Metadata-Version: 2.4
Name: dashai_stable_diffusion_v1_model_package
Version: 0.0.11
Summary: DashAI Stable Diffusion V1 Model Package for DashAI
Project-URL: Homepage, https://github.com/DashAISoftware/DashAI
Project-URL: Issues, https://github.com/DashAISoftware/DashAI/issues
Author: DashAI team
Author-email: dashaisoftware@gmail.com
Keywords: DashAI,Model
Requires-Python: >=3.8
Requires-Dist: diffusers>=0.20.0
Requires-Dist: torch>=2.0.1
Description-Content-Type: text/markdown

# Stable Diffusion V1 Model Plugin for dashAI

This plugin integrates CompVis' Stable Diffusion 1.x text-to-image models into the dashAI platform using the Hugging Face `diffusers` backend. It enables flexible image generation from text prompts, with configurable inference and sampling parameters.

## Included Models

### 1. Stable Diffusion v1-1
Original 1.x checkpoint from CompVis, based on [`CompVis/stable-diffusion-v1-1`](https://huggingface.co/CompVis/stable-diffusion-v1-1).

### 2. Stable Diffusion v1-2
Continued training over v1-1, based on [`CompVis/stable-diffusion-v1-2`](https://huggingface.co/CompVis/stable-diffusion-v1-2).

### 3. Stable Diffusion v1-3
Continued training with classifier-free guidance, based on [`CompVis/stable-diffusion-v1-3`](https://huggingface.co/CompVis/stable-diffusion-v1-3).

### 4. Stable Diffusion v1-4
Most refined 1.x checkpoint, based on [`CompVis/stable-diffusion-v1-4`](https://huggingface.co/CompVis/stable-diffusion-v1-4).

All models are latent text-to-image diffusion models, designed for high-quality image generation and compatible with CPU or GPU inference through `diffusers`.

## About Stable Diffusion

Stable Diffusion is a latent text-to-image diffusion model developed by CompVis, Stability AI, and Runway. It generates photo-realistic images from natural language descriptions.

Key features of Stable Diffusion 1.x models:

- Latent diffusion: generation runs in a compressed latent space for efficiency
- Text-conditioned: images are guided by natural language prompts
- Negative prompting: steer the model away from unwanted elements
- Open weights: freely available checkpoints
- Configurable resolution: supports a range of output dimensions

Stable Diffusion is designed for deployment on laptops, desktops, and cloud infrastructure, making advanced image generation more accessible.

## Features

- Image generation from text prompts
- Multiple 1.x checkpoints selectable at runtime
- Reproducible generation via fixed random seeds
- Configurable generation parameters:
  - `negative_prompt`: elements to avoid in the image
  - `num_inference_steps`: number of denoising steps
  - `guidance_scale`: how strongly the model follows the prompt
  - `width` / `height`: output image dimensions
  - `num_images_per_prompt`: number of images per prompt
  - `device`: detected GPU (e.g. `"GPU 0: NVIDIA RTX 3090 - Compute Capability 8.6"`) or `"CPU"`

## Model Parameters

| Parameter | Description | Default |
|-----------|-------------|---------|
| `model_name` | Stable Diffusion 1.x checkpoint from Hugging Face | `"CompVis/stable-diffusion-v1-1"` |
| `negative_prompt` | Text prompt for elements to avoid in the image | `""` (optional) |
| `num_inference_steps` | Number of denoising steps (higher = better quality, slower) | 15 |
| `guidance_scale` | How strongly the model follows the prompt | 3.5 |
| `seed` | Random seed for reproducibility (negative = random) | -1 |
| `width` | Width of the generated image (multiple of 8) | 512 |
| `height` | Height of the generated image (multiple of 8) | 512 |
| `num_images_per_prompt` | Number of images to generate per prompt | 1 |
| `device` | Inference device — a detected GPU by name (e.g. `"GPU 0: NVIDIA RTX 3090 - Compute Capability 8.6"`) or `"CPU"` | First detected GPU if available, else `"CPU"` |

## Requirements

- `diffusers`
- `torch`
- Model files from Hugging Face:
  - [`CompVis/stable-diffusion-v1-1`](https://huggingface.co/CompVis/stable-diffusion-v1-1)
  - [`CompVis/stable-diffusion-v1-2`](https://huggingface.co/CompVis/stable-diffusion-v1-2)
  - [`CompVis/stable-diffusion-v1-3`](https://huggingface.co/CompVis/stable-diffusion-v1-3)
  - [`CompVis/stable-diffusion-v1-4`](https://huggingface.co/CompVis/stable-diffusion-v1-4)

## Notes

This plugin uses the Hugging Face `diffusers` pipeline for text-to-image generation.

The models are pretrained for inference and are not designed for fine-tuning through this plugin.
Image dimensions must be multiples of 8.
